stable-diffusion-image-generation
Generates images with Stable Diffusion models (SD 1.5, SDXL, SD 3.0, Flux) through the HuggingFace Diffusers library, covering text-to-image, image-to-image, inpainting, outpainting, ControlNet conditioning, LoRA adapters, scheduler swapping, and GPU memory optimization. Use when generating images from text prompts, transforming or restyling an existing image, filling masked regions of an image, adding spatial control from edges, poses, or depth maps, loading LoRA style adapters, or fixing out-of-memory and black-image errors in a Diffusers pipeline. Not for API-only generation without a GPU, where DALL-E 3 fits better.
- Version
- 1.0.0
- License
- MIT
Pinned to revision df088027ff23, so it is the text this page describes rather than whatever the author pushed since.
Files
- skills/stable-diffusion-image-generation/SKILL.md
- skills/stable-diffusion-image-generation/references/advanced-usage.md
- skills/stable-diffusion-image-generation/references/controlnet.md
- skills/stable-diffusion-image-generation/references/generation-parameters.md
- skills/stable-diffusion-image-generation/references/troubleshooting.md
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